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README.md

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+# Stable-Diffusion-XL-Base-1.0
+
+## Model Description
+**Stable Diffusion XL Base 1.0 (SDXL 1.0)** is a foundation text-to-image model released by Stability AI.  
+It is the flagship successor to Stable Diffusion 2.1, designed for photorealism, artistic flexibility, and high-resolution generation.  
+
+SDXL 1.0 is a latent diffusion model trained on a broad dataset of images and captions. Compared to prior versions, it improves prompt alignment, visual coherence, and output quality, especially in complex scenes and detailed compositions.
+
+## Features
+- **High fidelity image generation**: sharper details and improved realism.
+- **Flexible style range**: from photorealistic renders to artistic illustration.
+- **Better prompt alignment**: improved understanding of nuanced or multi-concept prompts.
+- **High resolution support**: natively trained for 1024×1024 images.
+- **Compositional strength**: more accurate handling of multiple subjects and fine object placement.
+
+## Use Cases
+- Creative content generation (illustrations, art, concept design)
+- Product mockups and marketing visuals
+- Character and environment ideation
+- Storyboarding and visual storytelling
+- Research in generative imaging
+
+## Inputs and Outputs
+**Input**:
+- Text prompts (descriptions, concepts, artistic directions)
+- Optional negative prompts to avoid undesired elements
+
+**Output**:
+- Generated image(s) matching the prompt
+- Default resolution: 1024×1024 pixels
+
+## License
+- Licensed under: [CreativeML Open RAIL++-M License](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0/blob/main/LICENSE)
+
+## References
+- Model card: [https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0)